- Stable diffusion reddit It's extremely important for fine-tuning purposes and understanding the text-to-image space. But how much better? Asking as someone who wants to buy a gaming laptop (travelling so want something portable) with a video card (GPU or eGPU) to do some rendering, mostly to make large amounts of cartoons and generate idea starting points, train it partially on my own data, etc. Since a lot of people who are new to stable diffusion or other related projects struggle with finding the right prompts to get good results, I started a small cheat sheet with my personal templates to start. Stable Diffusion web UI: A browser interface based on Gradio library for Stable Diffusion runs locally (includes GFPGAN, Textual inversion, Mask painting, RealESRGAN and many more features) Easy Diffusion is a Stable Diffusion UI that is simple to install and easy to use with no hassle. Negative prompts for anatomy etc. Comparative between different styles using the same seed image at low noise levels /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. I completely disagree. AI models come in two types : pretrained, and fine-tunes. Even if you could somehow force Stable Diffusion to only use Super 8 film sources, you probably wouldn't want to. If you're using some web service, then very obviously that web host has access to the pics you generate and the prompts you enter, and may be doing something with said images. It's incredibly unoptimized, lacks any nice gui, etc. The two keys to getting what you want out of Stable Diffusion are to find the right seed, and to find the right prompt. 5) along with most community-made models. ckpt or . I know this is likely an overly often-asked question, but I find myself inspired to use Stable Diffusion, see all these fantastic posts of people using it, and try downloading it, but it never seems to work. This guide assumes that you are already familiar with Automatic111 interface and Stable Diffusion terminology, otherwise see this wiki page. If you are running stable diffusion on your local machine, your images are not going anywhere. How-to train a custom model and resources on doing so. Hi guys! I'm very interested in how to draw more than one person, for example 5 or more, that is, to choose exactly how many… Stable Diffusion Samplers: A Comprehensive Guide (stable-diffusion-art. 20-30 or so seems to generate a more complete looking image in a comic- digital painting style. Dec 23, 2024 · Users share their questions and experiences on how to fine tune Stable Diffusion XL, a generative AI model for images, using loras, dreambooths, or other methods. 1-0. Specially trained towards certain subjects and/or styles. Getting a single sample and using a lackluster prompt will almost always result in a terrible result, even with a lot of steps. 1. That is a bit of a massive number. " Some of this is easy to drop. Make sure that is high You should reread u/PlanetaryDecay comment. After a dozen hundred generations over the last week, I've come to treat Stable Diffusion a lot like working in a Darkroom, going from General to Specific. Simply choose the category you want, copy the prompt and update as needed. don't work. If you get only 1 great image in 100 then something is up with your prompts or the model you are using. I'm using Stable Diffusion locally and love it, but I'm also trying to figure out a method to do a complete offline install. In this article I have compiled ALL the optimizations available for Stable Diffusion XL (although most of them also work for other versions). We used Controlnet in Deforum to get similar results as Warpfusion or Batch Img2Img. One common way is to use the loss function. They also comment on the VRAM and speed requirements of Flux and other models. Technical details regarding Stable Diffusion samplers, confirmed by Katherine: - DDIM and PLMS are originally the Latent Diffusion repo DDIM was implemented by CompVis group and was default (slightly different update rule than the samplers below, eqn 15 in DDIM paper is the update rule vs solving eqn 14's ODE directly) TUTORIAL As someone who works in AI and does hobby game dev, Stable Diffusion has had such a huge impact on my workflow. 5 I generate in A1111 and complete any Inpainting or Outpainting, Then I use Comfy to upscale and face restore. But latent, by definition, means unobservable. That said, you're probably not going to want to run that. Seeds are crucial for understanding how Stable Diffusion interprets prompts and allow for controlled experimentation. I think I'm ready to upgrade to a better service, mostly for better resolutions, less wait times, and more options. For example, I might want to have a portrait I've taken of someone altered to make it look like a Picasso painting. If I have a lot of Loras in my prompt, I typically go with Euler because it’s better for the higher steps, and if I’m going for something a little less complex, I usually go for the Karras. 3" you can do the following: Writing (apple) puts more weight on the word apple. 4/5 has some flaws like hands, setting the resolution too big for initial creation etc. Base standard SD1. Aspect Ratios and CFG Scale: Aspect Ratios: The ratio of an image's width to its height, which has a significant impact on image generation. A1111 is another UI that requires you to know a few Git commands and some command line arguments but has a lot of community-created extensions that extend the usability quite a lot. "This is a guide for X Stable diffusion model that is currently released, and available here: link" and right there, that's three of the Ws. This is found under the "extras" tab in Automatic1111 Hope that makes sense (and answers your question). I'd recommend looking at other still film stocks of the same era, or possibly trying out using components of the style your trying to achieve like (light_leak:0. This is the workflow which gives me best results when trying to get a very specific end result. I've read it can work on 6gb of Nvidia VRAM, but works best on 12 or more gb. This means that the model is no longer changing significantly, and the generated images are becoming more realistic. Let's start with the first one. When inpainting, you can raise the resolution higher than the original image, and the results are more detailed. 5, that model that was released before the "safety stuff" was applied to it ? you may not care if models are used in bad ways but I can tell you it gave me sleepless nights. A sufficiently advanced model should be able to recognize the perspective of the viewport (maybe from the surrounding scene, which this model cuts out) and use that to not create incorrect perspectives, or simply recognize the object as a train and have the knowledge that trains don't look like that. Stable Diffusion not just works well on standard GPUs but also mining GPUs as well and it could be a cheaper alternative for those who are wanted a good or better GPU yet having much budget constraint for it. If you are really unlucky the out of memory will require a restart of the a1111 cmd window before it works with your new settings. I bought Nvidia P104 8GB GDDR5 GPU for $25 and it is fairly cheap for me. Now let's push past the random introductions, which aren't introductions A mix of Automatic1111 and ComfyUI. Here’s a data explorer for “Ghibli” images. I explain how they work and how to integrate them, compare the results and offer recommendations on which ones to use to get the most out of SDXL, as well as generate images with only 6 GB of graphics card memory. Merge Models. I'll try with different guidance scales to see what impact that makes. Install the Dynamic Thresholding extension So I was sitting here bored and had the idea of running some song lyrics to see what sort of pics I'd get, just for shits and gigs. Aug 12, 2024 · A user shares their thoughts on the misinformation and nuances of Flux, a new version of StableDiffusion, a text-to-image generation model. Some of this is "Duh you don't have to say it. So for example, if I have a 512x768 image, with a full body and smaller / zoomed out face, I inpaint the face, but change the res to 1024x1536, and it gives better detail and definition to the area I am for training models/lora the best guides are mostly on civitai under their articles section, also some on reddit, in fact sometimes just a random comment on reddit ends up giving a more interesting tip than any tutorial could, at least once you get the basics down, there are some tutorials on youtube but honestly I find these written guides DreamShaper: Best Stable Diffusion model for fantastical and illustration realms and sci-fi scenes. For this test we will review the impact that a seeds has on the overall color, and composition of an image, plus how to select a seed that will work best to conjure up the image you were envisioning. com) Choosing a sampler for Stable Diffusion (mccormickml. Seeds: Unique image IDs that help generate consistent images with slight variations. TLDR: Results 1, Results 2, Unprompted 1, Unprompted 2, links to checkpoints used at the bottom. My drawing skills are non-existant (as in I can't draw anything properly, not even write my name properly) and from what I've seen online, it is very possible to create a comic from Apr 3, 2024 · Learn how to write effective prompts for Stable Diffusion AI, a powerful image generation tool. On my part the artifacts were created by the tiles which are created with the tile upscales. While the synthetic (generated) captions were not used to train original SD models, they used the same CLIP models to check existing caption similarity and decide Right? On those expressions, small things really do make a world of a difference. 2) or (retro_color_grading:0. I am just getting started with video generation and any advice is appreciated. Comfy is great for VRAM intensive tasks including SDXL but it is a pain for Inpainting and outpainting. Hey, I love free stuff, use Stable Diffusion locally, but with that attitude the community is screwed long term. Most people use Automatic1111's webui which currently supports Stable Diffusion 1. The paper that gave one of the bases for modern upscaling was Super Resolution which proposed the SRGAN, used these indices for calculations. so which GUI in your opinion is the best (user friendly, has the most utilities, less buggy etc) personally, i am using… In the context of Stable Diffusion, converging means that the model is gradually approaching a stable state. However, sampling speed and memory constraints remain a major barrier to the practical adoption of diffusion models as the generation process for these models can be slow due to the need for iterative noise estimation using complex neural networks. For SD 1. The other post links ways to use Stable Diffusion locally on your own gpu. Use pre-trained Hypernetworks. I've been using several AI LLMs like vicuna, Stable Diffusion and training with a Radeon 6700 XT 12GB, in several Linux distributions (Fedora, Ubuntu, Arch) without any special driver installation, only installing ROCm with pip (python package installer). In other words, it is a black box where no one really knows what's exactly going on in that space. I feel like two 3090s at $600 or so apiece offers better value at 1200 than a $1600 4090, since I don't believe a 4090 actually comes close to being twice as fast. Here they are below. UI Plugins: Choose from a growing list of community-generated UI plugins, or write your own plugin to add features to the project! /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. ComfyUI is nice but it's a pain if you rely on inpainting a lot. Yes there is, there are 2 stats that can be used PSNR : Peak Signal to Noise Ratio and SSIM : Structural Similarity Index Measure. They suggested you open the new image in img2img(inpaint) and keep the original prompt, while setting the denoise strength really low, which probably means around 0. If we want a lot more innovation and investiment it will cost something. You can use your own list of styles, characters, objects or use the default ones which are already kinda huge. Hiya guy! ive collected around 100 or more really kinky and hot prompts which can work well and run all together with the script option in stable diffusion. It just sees a bag of pink or brown or whatever pixels. safetensors file, by placing it inside the models/stable-diffusion folder! Stable Diffusion 2. I'm talking - bring all required files on a Hard Drive to a laptop that has 0 connections, and making it work. Licensing: Stable Diffusion 3 Medium is open for personal and research use. 3. The best pony model I tested so far. r/StableDiffusion: /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper… Join the discussion and learn how to create AI art with Stable Diffusion, a powerful text-to-image generator. I'm a photographer and am interested in using Stable Diffusion to modify images I've made (rather than create new images from scratch). Hey, I'm currently developing a stable diffusion API as well. It’s great for people that have computers with weak gpus, don’t have computers, want a convenient way to use it, etc. 11 votes, 30 comments. A new Creator License enables professional users to utilize SD3 while supporting Stability's mission to democratize AI, maintaining a commitment to open AI. A finetune is a modification of an existing model. Today I discussed some new techniques on a livestream with a talented Deforum video maker. Folks developing can't eat air, and have to pay server costs. I really should have added my settings to the post. g. May I ask, what's your reasoning for not using the official API? 34 votes, 24 comments. 5 user for anime images and honestly was pretty wholly satisfied with it except for some few flaws like anatomy, taking forever to semi-correctly inpaint hands afterwards etc. . As noted in my test of seeds and clothing type, and again in my test of photography keywords, the choice you make in seed is almost as important as the words selected. So if there is no other way of reducing the blurs, then you should create a custom Target size for example 1072 x 1440 and the created Tile (under the upscaler you choose) should have the same size example 1072 x 1440, mask blur i had 8 and padding 42. Stable Diffusion is a latent Diffusion model involving latent space. Stable Video Diffusion, AnimateDiff, Lavie, Latte are some models I came across, but all are ~6 months old. See examples, links, and tips on creating realistic or artistic images. Always I get stuck at one step or another because I'm simply not all that tech savvy, despite having such an interest in these types of Although these images are quite small, the upscalers built into most versions of Stable Diffusion seem to do a good job of making your pictures bigger with options to smooth out flaws like wonky faces (use the GFPGAN or codeformer settings). It's not up yet, but you can feel free to join the discord for now It'll likely be cheaper than the official one. I'll do that now. It can be used entirely offline. Use custom VAE models. I update with new sites regularly and I do believe that my post is the largest collection of Stable Diffusion generation sites available. Why is it so challenging for Stable Diffusion 1. I use Stable Diffusion through Automatic1111's webui and a homemade software to navigate previous generations. 5). Find out the components, techniques, and tips for creating your desired image with keywords, art styles, and settings. I'm kinda new with Stable Diffusion and I'd like to find some good tutorial in order to achieve a specific goal: I'd like to be able to create a comic from scratch. Anything v5: Best Stable Diffusion model for anime styles and cartoonish appearance. Posted by u/AUTOMATIC1111 - 124 votes and 64 comments Here's my attempt to ELI5 how Stable Diffusion works: Billions of images are scraped from Pinterest, blogs, shopping portals, and other websites. Collect a dataset of images with the type of emotions you want, train a model, then you can blend it with another to get the emotions you want into the style you want. Nowadays, I almost exclusively use Krita AI Diffusion and feel pretty happy about it. After following these steps, you won't need to add "8K uhd highly detailed" to your prompts ever again: Install a photorealistic base model. To create a new piece of art, I generate a large number of intermediate images, which I then merge together and edit in Photoshop. Craft your prompt. x models (up to 1. Saldy my GPU has given up so if anyone is interested in using them. There are a few different ways to measure convergence in Stable Diffusion. Juggernaut XL: Best Stable Diffusion model for photography-style images/real photos. Additionally, I find it difficult to apply a style when the person in the image is an alien or some animal. That was interesting but I got curious about how well SD knew some of my old fave artists, and quickly realized that they (and I) are all a lot older now, so most of the pics are older folks, but occasionally it threw in some elements from the younger person, like Thanks for pointing this out. I feel like it's just a standard tool at the this point I wouldn't want to go without. Trial and error, probably errors until you find what works. I use cfg 3 and dpm sgm unifrom or euler(non acncestral) with 30 steps as sampler. Hi everyone, like many of you, I was amazed at the capabilities of Diffusion models, and it becomes a daily habit to generate a few images just for fun. 13 votes, 18 comments. It has REALLY good skin textures. This seems like Warpfusion, which has been the best method for getting stable (ha!) style transfer to videos with Stable Diffusion. This is the absolute most official, bare bones, basic code/model for Stable Diffusion. Learn how to improve your skills in using Stable Diffusion even if a beginner or expert. and of course Stable diffusion- automatic1111 build, but honestly I don't have loyalty to a specific build- I just like using what is currently the least buggiest and most convenient to use if that makes sense sorry for the late reply! Feel free to ask more questions if you have any! I wanted to share my latest exploration on Stable Diffusion - this time, image captioning. Really no problem my dude, just a copy paste and some irritability about everything having to be a damn video these days. This is my workflow others will do it better/different but it works Hi, neonsecret here I again spend the whole weekend creating a new UI for stable diffusion this one has all the features on one page, and I even made a video tutorial about how to use it. /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. Fast-forward a few weeks, and I've got you 475 artist-inspired styles, a little image dimension helper, a small list of art medium samples, and I just added an image metadata checker you can use offline and without starting Stable Diffusion. Find guides, tips, prompts, and links to online and offline tools for Stable Diffusion. 5 to apply the style consistently, even when using a Lora or a specific model? I've been experimenting with this for the past four months, and it works well with one model but not with another. I started messing with AI image stuff for the first time about a month ago, and after two days I realized the rate of learning/finding new things was just damn near exponential. I found it written in the example prompts of the stable diffusion pipeline used by the huggingface resource page and have used this style for my prompts ever since I do know that for some SD models, like "Realistic Vision 1. Krita's AI plugin is a perfect solution for me, which combines Krita's powerful layer-based image editing capabilities with the flexibility of ComfyUI. The problem is that it doesn't know what hands and other things are. In more mathematical terms, the process in latent space cannot be described by a function q(x). com) Can anyone explain differences between sampling methods and their uses […] ? (reddit) Can anyone offer a little guidance on the different Samplers? (reddit) What are all the different samplers (github. , HTML alt-text tags) and other fields. 1 support. true. com) Hereby, I present to you a Stable Diffusion Prompt Generator, which can randomly (or less randomly, depending on your inputs) build tons of prompts for you : Stable Diffusion Random Prompts Generator. Hey folks – I've put together a guide with all the learnings from the last week of experimentation with the SD3 model. - Setup - All images were generated with the following settings: Steps: 20 Sampler: DPM++ 2M Karras When you use Stable Diffusion, you use models, also called checkpoints. I created this for myself since I saw everyone using artists in prompts I didn't know and wanted to see what influence these names have. Custom Models: Use your own . Since a big Another trick I haven't seen mentioned, that I personally use. Go to “Make Animation”, upload your stable diffusion image through your photo library Pick a Video (or make one using the camera) to drive the animation Wait (or send a few more while you wait) You’ve got an animation Workflow option 2 (run it yourself): Now i know people say there isn't a master list of prompts that will get you magically get you perfect results and i know that and thats not quite what im looking for but i simply need help with prompts since im not really that descriptive especially when it comes to hairstyles and poses. Trains never look like that. Making a pretrained model is extremely expensive (you need multiple GPUs running full time for days), which is why research leaned towards finetunes. Try zonkey. Hey all, I've been really getting into Stable Diffusion lately but since I don't have the hardware I'm using free online sites. In general, for 99% of all the new fancy open source AI-stuff searching for "nameofthingyouwant github" on any search engine mostly takes you directly to the project where most of the time there's an official installation guide or some sort of explanation on how to use it. Hopefully some of you will find it useful. Sometimes when you come up with new models, it's necessary to go back to step 0 of the model and build it better from the ground up, which is what the people building Stable /r/StableDiffusion is back open after the protest of Reddit killing open API access, which will bankrupt app developers, hamper moderation, and exclude blind users from the site. Personally I will probably try to make a substantially larger emotions wordlist to see what else is out there that works, but in the end the best thing would be to have models trained from the ground up around the concept of character creation. Tokens trained on specific subjects and/or styles. 30-50 will be better for more detailed Edit: To be clear, the most important setting for good blending is the pixel/mask padding. We plan to incorporate other features like dreambooth and toggling of banned words in the future. What about Stable Diffusion Model 1. Get new upscalers here put them in your stable-diffusion-webui\models\ESRGAN folder. I was a big 1. When all else fails, the answer is “finetune your own model”. With SDXL picking up steam, I downloaded a swath of the most popular stable diffusion models on CivitAI to use for comparison against each other. Stable Diffusion and Disco Diffusion are both Diffusion models, but with significant changes on the learning/generative process, leading to dramatically different results. From what I've gathered from a less under the hood perspective: steps are a measure of how long you want the ai to work on an image (1 step would produce a image of noise while 10 might give you something starting to resemble an image but blurry/smudges/static. I’m a DPM 2M++ Karras at 25-30 steps or a Euler A 35-40 steps kind of guy depending on what I’m going for. Abstract Diffusion models have recently achieved great success in synthesizing diverse and high-fidelity images. These images are saved in a database along with their text descriptions (e. gsp gpcxt abvau aafaai lzteo rpznncfp kmppvhx eadr iwdiie wci whsa lpkpq cgrz rbezebps wdarxcq