The latest version of this model is Stable Diffusion XL, which has a larger UNet backbone network and can generate even higher quality images. Actually, It helps the generator understand what to avoid while creating the image. E. py --model stable_diffusion_xl --prompt "butterfly in new york, 4k, realistic" --clients 2 –requests 5 Examine the server log and metrics. to get started. You can harness these prompts as provided or adapt them as inspiration for your personalized phantom. This is where you'll find your prompt and negative prompt. As we all know, StabilityAI claims that the model is optimized for generating images from concise prompts. We deployed our model here; You can even run your trained model on segmind instead of the API; Fine-tuned model deployed as an API on Segmind. 2. I batch up the prompt and leave my PC for 30-40mins while it churns out the images. It is a parameter that tells the Stable Diffusion model what not to include in the generated image. Conclusion The training script is also similar to the Text-to-image training guide, but it’s been modified to support SDXL training. Here are the best prompts for Stable Diffusion XL collected from the community on Reddit and Discord: 📷 Stable Diffusion XL (SDXL) is an open-source diffusion model, the long waited upgrade to Stable Diffusion v2. Compared to Stable Diffusion V1 and V2, Stable Diffusion XL has made the following optimizations: Improvements have been made to the U-Net, VAE, and CLIP Text Encoder components of Stable Diffusion. Renowned for its versatility, this model is adept in many scenarios, offering impressive performance across diverse applications. In painting, the posture of a character is crucial for vividly expressing their life situation, emotional state, and mood. 0 model. In all seriousness, models can be trained using different data sets to excel in specific types of tasks. Tags vs Sentence. Jun 22, 2023 · Want to learn prompting techniques within Stable Diffusion to produce amazing results from your ideas? Well, look no further than this short, straight to the Oct 22, 2023 · For the purpose of the demonstration we will use the same prompt and use the XYZ Plot option in Stable Diffusion WebUI to plot both the camera framing and angle. Il est conçu pour servir de tutoriel, avec de nombreux exemples pour illustrer l’utilité ou le fonctionnement d’un paramètre. The architecture of the SDXL 1. or. 0. ago. a masterpiece of an ocean with waves and palm trees, sunset in We present SDXL, a latent diffusion model for text-to-image synthesis. Steps: 4-6. And with the built-in styles, it’s much easier to control the output. 0 specifically for pixel art. Also, for all the prompts below, I’ve purely used the SDXL 1. 0 as its base checkpoint, which is the cutting edge Stable Diffusion model at the time of writing, it is a versatile LoRA that can do both app Dec 27, 2022 · Well, you need to specify that. Written by: Milan K. Stable Diffusion XL (SDXL) is a powerful text-to-image generation model that iterates on the previous Stable Diffusion models in three key ways: the UNet is 3x larger and SDXL combines a second text encoder (OpenCLIP ViT-bigG/14) with the original text encoder to significantly increase the number of parameters. Now, let me start showing you the examples I wrote for you today. Start by loading up your Stable Diffusion interface (for AUTOMATIC1111, this is “user-web-ui. In this tutorial we will be going through the prompt formula for Stable Diffusion XL TURBO model and different parameter settings that will help you generate In this report, we discuss the art and science of prompt engineering for diffusion models primarily focussing on the Stable Diffusion family of models. Stable Diffusion XL is the latest iteration of the popular text-to-image generation model, offering impressive results. Unlikely that your image will look like their image. You can find a detailed guide on integrating our API's here: Stable Diffusion XL 1. It is created by Stability AI. Stable Diffusion XL output image can be improved by making use of a refiner as shown Prompt techniques. bat”). It provides a user-friendly way to interact with Stable Diffusion, an open-source text-to-image generation model. This applies to anything you want Stable Diffusion to produce, including landscapes. 1 comment. I wanted to change the style for two prompts. 3D rendering. Nov 20, 2023 · Best Stable Diffusion XL Photorealistic Prompts. It’s significantly better than previous Stable Diffusion models at realism. Find recommended settings, samplers, VAE, and 33 examples of prompt templates for various genres. ; prompt_2 (str or List[str], optional) — The prompt or prompts to be sent to the tokenizer_2 and text_encoder_2. prompt #1: character covered in a lush, living tapestry of flora, with vines and leaves that appear as a natural, enchanted armor. 8: Look at other people's prompts. selective focus, miniature effect, blurred background, highly detailed, vibrant, perspective control. Mar 7, 2024 · python3 client. Aug 25, 2023 · No longer must users resort to terms like “masterpiece” to yield high-quality results. Text prompts are used as the primary conditioning factor. Moreover, I Jul 2, 2023 · A good Stable Diffusion prompt should be: Clear and specific: Describe the subject and scene in detail to help the AI model generate accurate images. Table of Contents. 0” is a LoRA that is fine tuned to create app icons sized for iOS and Android apps. But you can also set its parameters to 768×768 or 512×512. We can even pass different parts of the same prompt to the text encoders. If not defined, one has to pass prompt_embeds. It's a versatile model that can generate diverse Mar 19, 2024 · Tips for good prompts. - the UI, model, image dimensions, seed and other factors determine if your image is going to look like their image. Example 1: An image of a beautiful sunset at the beach might be tagged by CLIP as: ocean, beach, sunset, waves, sand, palm trees, golden sky. Blog post about Stable Diffusion: In-detail blog post explaining Stable Diffusion. You have to be specific. Stable Diffusion is a Aug 6, 2023 · Learn how to use SDXL v1. Aug 19, 2023 · Ce guide a pour vocation de vous aider à maîtriser l'interface graphique d'AUTOMATIC1111. The Stable Diffusion XL (SDXL) model effectively comprises two distinct models working in tandem: ‍. It has a base resolution of 1024x1024 pixels. The jump from 768×768 in SD 2. The best prompts are detailed, specific, and well-structured to help the model realize your vision. Try out this prompt on OpenArt Model: Stable Diffusion Seed: 70455 Scale: 13 Steps: 25 Resolution: 512 x 512 Sampler: DDIM 8 Stable Diffusion XL works especially well with images between 768 and 1024. General info on Stable Diffusion - Info on other tasks that are powered by Stable . Switch between documentation themes. I’ve categorized the prompts into different categories since digital illustrations have various styles and forms. So, it’s like giving a little Learn how to create and use prompts for Stable Diffusion, ChatGPT and Midjourney models with PromptHero, the ultimate prompt search engine. We would like to show you a description here but the site won’t allow us. Deforum is a popular way to make a video from a text prompt. Jan 31, 2024 · Stable Diffusion Illustration Prompts. As a ballpark, most samplers should use around 20 to 40 steps for the best balance between quality and speed. 9 at the end. It uses a larger base model, and an additional refiner model to increase the quality of the base model’s output. Moving into detailed subject and scene description, the focus is on precision. links and info on use https:// rentry. You can use the free AI image generator on Stable Diffusion Online or search over 9 million Stable Diffusion prompts on Prompt Database. Conversely, a word inside square brackets receives reduced prominence. 0 model without any LORA models. The more vivid your mental image, the more detailed your prompt can be. Consider aspects such as the subject matter, setting, mood, color scheme, and lighting. You can keep adding descriptions of what you want, including accessorizing the cats in the pictures. ). Fooocus is a rethinking of Stable Diffusion and Midjourney’s designs: Learned from Stable Diffusion, the software is offline, open source, and free. What is Img2Img in Stable Diffusion. Stable Diffusion XL can pass a different prompt for each of the text encoders it was trained on as shown below. Feb 29, 2024 · Andrew. The Web UI offers various features, including generating images from text prompts (txt2img), image-to-image processing (img2img Oct 28, 2022 · Prompt Engineering Intro This is an example of a prompt and all the parameters Prompt: Funko pop superman figurine, made of plastic, product studio shot, on a white background, diffused lighting, centered. They are limited by the rather superficial knowledge of SD, but can probably give you a Rule 2. Might affect coherence and quality. This guide will focus on the code that is unique to the SDXL training script. Feb 12, 2024 · Here is our list of the best portrait prompts for Stable Diffusion: S. Conclusion . Aug 10, 2023 · SDXL 是 Stable Diffusion XL 的簡稱，顧名思義它的模型更肥一些，但相對的繪圖的能力也更好。 Stable Diffusion XL - SDXL 1. FlashAttention: XFormers flash attention can optimize your model even further with more speed and memory improvements. To access SDXL using Clipdrop, follow the steps below: Navigate to the official Stable Diffusion XL page on Clipdrop. Minimalist background, 50mm lens for focused facial expression. In order to increase emphasis on a word or phrase, add a + or number between 1. Compared to previous versions of Stable Diffusion, SDXL leverages a three times larger UNet backbone: The increase of model parameters is mainly due to more attention blocks and a larger cross-attention context as SDXL uses a second text encoder. It's designed for designers, artists, and creatives who need quick and easy image creation. Keep reading to start creating. 1st row prompt: <framing>, BREAK a gorgeous woman with short dark hair, new york city, <framing>, sunset Prompts On the first tab, txt2img, you'll be making most of your images. Stable Diffusion is not like Midjourney or other popular image generation software, you can't just ask it what you want. Here, the use of text weights in prompts becomes important, allowing for emphasis on certain elements within the scene. Fashion Model. Prompt: Monochromatic portrait of a model with dramatic makeup, chiaroscuro lighting for a classic art feel. Relevant: Use relevant keywords and phrases that are related to the subject and Sep 23, 2023 · tilt-shift photo of {prompt} . Jan 22, 2024 · Approach 1: SD XL Generated App Icons. These prompts guide the image generation process, ensuring that the produced images align with the text's description. No Account Required! Stable Diffusion Online is a free Artificial Intelligence image generator that efficiently creates high-quality images from simple text prompts. Good balance between creativity and guided generation. A Detailed Stable Diffusion Pose Prompt Guide with 15 Examples. Dreambooth - Quickly customize the model by fine-tuning it. Whether you're looking to visualize concepts, explore new creative avenues, or enhance Artist Inspired Styles. I’ve covered vector art prompts, pencil illustration prompts, 3D illustration prompts, cartoon prompts, caricature prompts, fantasy illustration prompts, retro illustration prompts, and my favorite, isometric illustration prompts in this Mar 8, 2024 · This guide elucidates prime prompt tips and embodies 15 unique SDXL prompts that span a broad stylistic spectrum. She wears a medieval dress. Stable Diffusion image 1 using 3D rendering. Oct 31, 2023 · A negative prompt for SDXL is like giving it a description of what you don’t want to see in the picture. One of the standout features of this model is its ability to create prompts based on a keyword. Negative prompting influences the generation process by acting as a high-dimension anchor, which Feb 20, 2024 · Stable Diffusion Prompts Examples. A separate Refiner model based on Latent has been Apr 3, 2024 · Here in our prompt, I used “3D Rendering” as my medium. Running Stable Diffusion Locally. Par exemple, l'ajout de "ugly" (laid), "disformed” (difforme) comme negative prompt peut aider à obtenir des images plus These tips and prompting styles will work with any model that directly uses pony diffusion v6 xl, like autismix pony for example. In Stable Diffusion, however, triple parentheses amplify a word’s The following prompts are supposed to give an easier entry into getting good results in using Stable Diffusion. 0 to 1024×1024 in SDXL represents a significant increase in the number of pixels – nearly Prompt examples : Prompt: cartoon character of a person with a hoodie , in style of cytus and deemo, ork, gold chains, realistic anime cat, dripping black goo, lineage revolution style, thug life, cute anthropomorphic bunny, balrog, arknights, aliased, very buff, black and red and yellow paint, painting illustration collage style, character Oct 30, 2023 · Today, after Stable Diffusion XL is out, the model understands prompts much better. Stable Diffusion XL. The higher resolution enables far greater detail and clarity in generated imagery. 0_0. Aug 4, 2023 · Undoubtedly, the emergence of Stable Diffusion XL has marked a milestone in the history of natural language processing and image generation, taking us a step closer to something that already scares… Sep 16, 2023 · In this comprehensive guide, we’ll walk you through setting up the software, using the color sketch tool, and leveraging Img2Img to turn amateur sketches into professional artwork. Related: Best Stable Diffusion Anime Prompts. Feb 22, 2024 · Stable Diffusion XL 1. co/ponyxl_loras_n_stuff was used to source as well as the purplesmart. 1-768. In this blog, we'll explore captivating prompts, each unlocking the potential of the SDXL 1. . Through posture, artists can profoundly depict a character's personality, emotions, and inner world, providing viewers with a deeper understanding of the artwork. Mar 29, 2024 · The fundamental usage of Stable Diffusion models is in the text-to-image format. The “ Icons. Jul 6, 2024 · You should see two nodes labeled CLIP Text Encode (Prompt). This visualization forms the foundation of your prompt. Read the SDXL guide for a more detailed walkthrough of how to use this model, and other techniques it uses to produce high quality images. You can use them as is or as a starting point to create your own. Stable Diffusion Portrait Prompts. 9: Good luck, and always be testing! Jan 25, 2024 · Poised, mid-dance movement, creating a sense of fluidity and grace. Since it is open source and anyone who has 5GB of GPU VRAM can download it (and Emad Mostaque Step 1. The structure of the prompt Parameters . Sign Up. Aug 5, 2023 · Stable Diffusion XL can produce images at a resolution of up to 1024×1024 pixels, compared to 512×512 for SD 1. 3. This model allows for image variations and mixing operations as described in Hierarchical Text-Conditional Image Generation with CLIP Latents, and, thanks to its modularity, can be combined with other models such as KARLO. Sep 16, 2023 · A negative prompt is a way to use Stable Diffusion in a way that allows the user to specify what he doesn’t want to see, without any extra input. photograph should not be used with van Gogh. It might take a few minutes to load the model fully. Prompts are important because they describe what you want a diffusion model to generate. Apocalyptic Space Station. May 16, 2024 · Recommended Setting Hyper Version (VAE is baked in): Res: 832*1216 (Any SDXL Res will work fine) Sampler: DPM++ SDE Karras. Nov 23, 2023 · I really like the fact that there is a preset style in Stable Diffusion XL 1. 1. Diffusion models are taught by introducing additional pixels called noise into the image data. Nov 12, 2023 · Stable Diffusion Prompts for NFT Art (14 Prompt Examples) In this guide I'll show you how to use a generative AI model like Stable Diffusion XL to create NFT art and characters. If you're unsure how to select this style, I'll show you in the image below. First, we establish a simple image generation workflow using the 🧨 Diffusers library by 🤗 HuggingFace. In every step, the U-net in Stable Diffusion will use the prompt to guide the refinement of noise into a picture. That's why I wrote two prompts to go along with that specific style. GBJI • 1 min. It’s worth mentioning that previous Fooocus. In order to reduce emphasis on a word or phrase, add a - or number between 0 and 0. Photo of a man with a mustache and a suit, plain background, portrait style. Here, for each type of photorealistic image, I will first provide a brief explanation of that type. Not Found. Stable Diffusion separates the imaging process into a diffusion process at runtime. However, a great prompt can go a long way in generating the best output. Depending on the hardware available to you, this can be very computationally intensive and it may not run on a Aug 30, 2022 · In this guide, you will learn how to write prompts by example. 0 Model - Stable Diffusion XL Stable Diffusion XL or SDXL is the latest image generation model that is tailored towards more photorealistic outputs… Sep 8, 2023 · In this article, you will find some prompt tips and 15 SDXL prompts with different styles. 9vae. It starts by creating functions to tokenize the prompts to calculate the prompt embeddings, and to compute the image embeddings with the VAE. CFG 7 – 10: Recommended for most prompts. Setup Python and Pip. The following prompts are mostly collected from different discord servers, websites, fabricated and then modified Jan 12, 2023 · Stable Diffusion is a text-to-image model that uses a frozen CLIP ViT-L/14 text encoder. For example, if you’re asking for a picture of a happy dog, you should use a negative prompt, like “No sad dogs”. 05 times. The default weight = 1. The prompt affects the output for a trivial reason. Adding noise in a specific order governed by Gaussian distribution concepts is Oct 29, 2023 · This Stable Diffusion XL or SDXL prompt guide aims to provide a comprehensive understanding of various aspects related to the SDXL model from prompt anatomy to styles. In Stable Diffusion, models determine what art style they can Dec 24, 2023 · Stable Diffusion XL (SDXL) is a powerful text-to-image generation model. CFG 16 – 20: Not generally recommended unless the prompt is well-detailed. 5 and 768×768 for SD 2. Prompt weighting basics. 1, Hugging Face) at 768x768 resolution, based on SD2. SDXL’s UNet is 3x larger and the model adds a second text encoder to the architecture. May 27, 2023 · En outre, les prompts négatifs sont particulièrement importants dans les modèles Stable Diffusion v2. Now Stable Diffusion returns all grey cats. Nov 30, 2023 · prompt #10: enchanting and colorful wonderland with fantastical creatures, whimsical landscapes, and playful adventures. This compendium, which distills insights gleaned from a multitude of experiments and the collective wisdom of fellow Stable Diffusion aficionados, endeavors to be a In Stable Diffusion, wrapping a word with triple parentheses ( ( (word))) boosts its weight by 1. ← ControlNet with Stable Diffusion 3 ControlNet-XS →. No. Copy it to your favorite word processor, and apply it the same way as before, by pasting it into the Prompt field and clicking the blue arrow button under Generate. Intense, contemplative look, blending modern fashion with classic art aesthetics. Concise: Use concise language and avoid unnecessary words that may confuse the model or dilute the intended meaning. If you’re eager to dive into the world of AI-generated art using Stable Diffusion XL, this guide will help you get started. Jan 29, 2024 · Food Photography. It saves you time and is great for quickly fixing common issues like garbled faces. instead. It wasn't long before many celebrities and Jul 22, 2023 · After Detailer (adetailer) is a Stable Diffusion Automatic11111 web-UI extension that automates inpainting and more. fashion editorial, a female model with blonde hair, wearing a colorful dress. We also demonstrate the use of the Weights & Biases integration for Diffusers to Stable Diffusion models. Feb 24, 2024 · Here’s my list of the best SDXL prompts. However, you can use Lora models to generate images with a specific style, object, or setting. 5B parameter base model and a 6. We’re on a journey to advance and democratize artificial intelligence through open source and open science. Add a Comment. Embarking on a journey with Stable Diffusion prompts necessitates an exploratory approach towards crafting veritably articulate and distinctly specified prompts. Very specific. The second way is to stylize a video using Stable Diffusion. Software. L'ajout de quelques prompts négatifs peut améliorer considérablement l'esthétique des images générées. Aug 3, 2023 · The Stable Diffusion model SDXL 1. In this post, you will learn how it works, how to use it, and some common use cases. Initially, the base model is deployed to establish the overall composition of the image. Img2Img is a cutting-edge technique that generates new images from an input image and a corresponding text prompt. The model distinguishes between concepts like “The Red Square” (a renowned location) and a “red square” (a geometric shape). Prompt: A beautiful ((Ukrainian Girl)) with very long straight hair, full lips, a gentle look, and very light white skin. Please post them in plain text over here with their accompanying pictures - that would be much more useful and reach more people. Faster examples with accelerated inference. NFTs seemingly came out of nowhere and skyrocketed in popularity. You have probably seen one of them on social media. N'hésitez pas à le bookmarquer pour le consulter également comme un manuel de référence. To begin, envision the image you wish to create. Clipdrop is a fantastic way to test out Stable Diffusion XL, for free! They do offer a pro version for faster image generation, but even on the free version the images don’t take very long at all to be generated. Essentially, when a word is enclosed in parentheses, the model emphasizes it more in its output. Examples of prompts for the Stable Diffusion process. Stable UnCLIP 2. 5 Upscale. Snowy Monster. Note that I use the Clipdrop platform by Stability AI for access to Stable Diffusion XL. I think using a video to present prompts is just about the worst medium you could use. Asian Zen Garden. Stable Diffusion XL (SDXL) SDXL is a more powerful version of the Stable Diffusion model. It looks like this. Redmond – App Icons Lora for SD XL 1. In Stable Diffusion v1, most users do that with multiple custom models. The artist’s name is a very strong style modifier. CFG 10 – 15: When you’re sure that your prompt is detailed and very clear on what you want the image to look like. While having an overview is helpful, keep in mind that these styles only imitate certain aspects of the artist's work (color, medium, location, etc. Simple prompts can already lead to good outcomes, but sometimes it's in the details on what makes an image believable. Be detailed and specific when describing the subject. Use Detailed Subjects and Scenes to Make Your Stable Diffusion Prompts More Specific. 1. 0) Image folder: <path to your image folder> Output folder: <path to the Cool Stable Diffusion XL prompts. In this list, you’ll find various styles you can try with SDXL models. 0 is Stable Diffusion's next-generation model. Fooocus is an image generating software (based on Gradio ). This will let you run the model from your PC. prompt (str or List[str], optional) — The prompt or prompts to guide the image generation. 35 Denoise + 1. Use an appropriate medium type consistent with the artist. But if you need to discover more image styles, you can check out this list where I covered 80+ Stable Diffusion styles. - keep a file of prompt ideas that you have copied and try them out. You need to know that the model is a switchable part of AI where magic is stored. 0 prompts and negative prompts to create stunning images with different styles. You can use the syntax (keyword:weight) to control the weight of the keyword. Sep 22, 2023 · Option 1: Every time you generate an image, this text block is generated below your image. If you want to run Stable Diffusion locally, you can follow these simple steps. On the checkpoint tab in the top-left, select the new “sd_xl_base” checkpoint/model. Use multiple brackets () to increase its strength and [] to reduce. Oct 24, 2023 · There is an option in Stable Diffusion that enables you to select this aspect ratio. It reminds me the moment when Stability AI Stable Diffusion XL (SDXL) is a larger and more powerful iteration of the Stable Diffusion model, capable of producing higher resolution images. 0 can be accessed and used at no cost. Feb 11, 2024 · Folders and source model Source model: sd_xl_base_1. Dec 18, 2023 · Enter any default prompt and click Import, the model will be imported in 2-3 minutes. Stable Diffusion XL enables us to create gorgeous images with shorter descriptive prompts, as well as generate words within images. Jan 30, 2024 · Stable Diffusion is a text-to-image model, powered by AI, that uses deep learning to generate high-quality images from text. Learned from Midjourney, the manual tweaking is not needed, and users only need to focus on the prompts and images. Stable Diffusion XL (SDXL) 1. 1 and 2 at the end. HiRes: 4xNMKD-Siax_200k with 3 Steps and 0. I ultimately decided that I was going to use the 'Fantasy Art' preset style in Stable Diffusion XL. Stable Diffusion is similarly powerful to DALL-E 2, but open source, and open to the public through Dream Studio, where anyone gets 50 free uses just by signing up with an email address. 0 model is built on innovation, comprising a 3. Best. Step 2. 6B parameter Jan 27, 2024 · Is Lora required to create landscape art in Stable Diffusion? Lora models aren’t required to create landscape art. So, that’s our list of the best landscape prompts for Stable Diffusion. Similarly, with Invoke AI, you just select the new sdxl model. Option 2: Install the extension stable-diffusion-webui-state. ai discord score_9, score_8_up, score_7_up, score_6_up, score_5_up, score_4_up, just describe what you want Collaborate on models, datasets and Spaces. May 5, 2024 · However, the effect of step count depends on the sampler chosen. But crafting a great prompt takes time and effort and sometimes it may not be enough because language and words can be imprecise. The more weight you add, the greater the risk of lower quality there is. 0 API Guide. Stable Diffusion image 2 using 3D rendering. Mar 16, 2024 · For local AI image generation, it’s hard to beat Stable Diffusion. 0 is the latest model in the Stable Diffusion family of text-to-image models from Stability AI. Enter your prompt in the top one and your negative prompt in the bottom one. g. Use "Cute grey cats" as your prompt instead. I generate images at 1024×1024 and resize them to speed up my webpage. Summary. I'll be using that style for every prompt featured in this article. CFG: 1-2 (recommend 2 for a bit negative prompt affection) Negative: Only working slightly on CFG 2. It can create images in variety of aspect ratios without any problems. Since it uses SD XL 1. Begin by confirming the installation of Python, as Python is the backbone of this project then, install or update Pip, the essential package manager for Python, to ensure all dependencies for Stable Diffusion XL can be managed effectively. With dynamic batching enabled, concurrent requests, and info-level logging, the example prints out additional information about the number of prompts included in each request to the TensorRT CLIP can be used to create detailed prompts for stable diffusion models by providing relevant tags that describe the content of an image. I mentioned the Stable Diffusion XL model a few times in this guide. We advise a fragmented Prompt such as: “ 1girl, schoolgirl, white uniform “ rather than a full sentence like: “ a schoolgirl in white uniform “ Even though they give very similar results with very small prompts, long sentence-type prompts are prone to be partially dismissed or get interrupted by unintended filler words. Insights on Prompting SDXL Embarking on the SDXL model journey requires an understanding of its intrinsic nuances. Stable Diffusion Web UI is a browser interface based on the Gradio library for Stable Diffusion. prompt: “📸 Portrait of an aged Asian warrior chief 🌟, tribal panther makeup 🐾, side profile, intense gaze 👀, 50mm portrait photography 📷, dramatic rim lighting 🌅 –beta –ar 2:3 –beta –upbeta –upbeta”. blurry, noisy, deformed, flat, low contrast, unrealistic, oversaturated, underexposed. Dec 21, 2023 · Step-by-Step Guide for Stable Diffusion XL Setup. Mar 8, 2024 · This guide elucidates prime prompt tips and embodies 15 unique SDXL prompts that span a broad stylistic spectrum. New stable diffusion finetune ( Stable unCLIP 2. The UNext is 3x larger. Golden Retriever. 500. safetensors (you can also use stable-diffusion-xl-base-1. Following this, an optional refiner model can be applied, which is responsible for adding more intricate details to the image. Published on: November 12, 2023. Let’s make some images Mar 19, 2024 · There are two main ways to make videos with Stable Diffusion: (1) from a text prompt and (2) from another video. The CLIP Text Enode node first converts the prompt into tokens and then encodes them into embeddings with the text encoder. A collection of what Stable Diffusion imagines these artists' styles look like. sc vq ef in oy ai sp zw jr nv